README.md

---
license: openrail++
language:
- en
tags:
  - text-to-image
  - stable-diffusion
  - lora
  - safetensors
  - diffusers
  - stable-diffusion-xl
base_model: Linaqruf/animagine-xl
inference: 
  parameter:
    negative_prompt: lowres, bad anatomy, bad hands, text, error, missing fingers, extra digit, fewer digits, cropped, worst quality, low quality, normal quality, jpeg artifacts, signature, watermark, username, blurry
widget:
- text: >-
    face focus, cute, masterpiece, best quality, 1girl, green hair, sweater, looking at viewer, upper body, beanie, outdoors, night, turtleneck
  example_title: example 1girl
- text: >-
    face focus, bishounen, masterpiece, best quality, 1boy, green hair, sweater, looking at viewer, upper body, beanie, outdoors, night, turtleneck
  example_title: example 1boy
datasets:
- Linaqruf/sdxl-dataset
---

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<h1 class="title"><span>Pastel Anime LoRA for SDXL</span></h1>
<h2 class="subtitle"><span>TRAINED WITH </span><a href="https://huggingface.co/Linaqruf/animagine-xl"><span>ANIMAGINE XL</span></a></h2>

<hr>

<table class="custom-table">
  <tr>
    <td>
      <a href="https://huggingface.co/Linaqruf/pastel-anime-xl-lora/blob/main/samples/xl_output_upscaled_00001_.png">
        <img class="custom-image" src="https://huggingface.co/Linaqruf/pastel-anime-xl-lora/resolve/main/samples/xl_output_upscaled_00001_.png" alt="sample1">
      </a>
    </td>
    <td>
      <a href="https://huggingface.co/Linaqruf/pastel-anime-xl-lora/blob/main/samples/xl_output_upscaled_00006_.png">
        <img class="custom-image" src="https://huggingface.co/Linaqruf/pastel-anime-xl-lora/resolve/main/samples/xl_output_upscaled_00006_.png" alt="sample2">
      </a>
    </td>
  </tr>
</table>

<hr>

## Overview

**Pastel Anime LoRA for SDXL** is a high-resolution, Low-Rank Adaptation model for Stable Diffusion XL. The model has been fine-tuned using a learning rate of 1e-5 over 1300 global steps with a batch size of 24 on a curated dataset of superior-quality anime-style images. This model is derived from Animagine XL.

Like other anime-style Stable Diffusion models, it also supports Danbooru tags to generate images.

e.g. _**face focus, cute, masterpiece, best quality, 1girl, green hair, sweater, looking at viewer, upper body, beanie, outdoors, night, turtleneck**_


<hr>

## Model Details

- **Developed by:** [Linaqruf](https://github.com/Linaqruf)
- **Model type:** Low-rank adaptation of diffusion-based text-to-image generative model
- **Model Description:** This is a small model that should be used with big model and can be used to generate and modify high quality anime-themed images based on text prompts. 
- **License:** [CreativeML Open RAIL++-M License](https://huggingface.co/stabilityai/stable-diffusion-2/blob/main/LICENSE-MODEL)
- **Finetuned from model:** [Animagine XL](https://huggingface.co/Linaqruf/animagine-xl)

<hr>

## 🧨 Diffusers 

Make sure to upgrade diffusers to >= 0.18.2:
```
pip install diffusers --upgrade
```

In addition make sure to install `transformers`, `safetensors`, `accelerate` as well as the invisible watermark:
```
pip install invisible_watermark transformers accelerate safetensors
```

Running the pipeline (if you don't swap the scheduler it will run with the default **EulerDiscreteScheduler** in this example we are swapping it to **EulerAncestralDiscreteScheduler**:
```py
import torch
from torch import autocast
from diffusers import StableDiffusionXLPipeline, EulerAncestralDiscreteScheduler

base_model = "Linaqruf/animagine-xl"
lora_model_id = "Linaqruf/pastel-anime-xl-lora"
lora_filename = "pastel-anime-xl.safetensors"

pipe = StableDiffusionXLPipeline.from_pretrained(
    model, 
    torch_dtype=torch.float16, 
    use_safetensors=True, 
    variant="fp16"
    )

pipe.scheduler = EulerAncestralDiscreteScheduler.from_config(pipe.scheduler.config)
pipe.to('cuda')

pipe.load_lora_weights(lora_model_id, weight_name=lora_filename)

prompt = "face focus, cute, masterpiece, best quality, 1girl, green hair, sweater, looking at viewer, upper body, beanie, outdoors, night, turtleneck"
negative_prompt = "lowres, bad anatomy, bad hands, text, error, missing fingers, extra digit, fewer digits, cropped, worst quality, low quality, normal quality, jpeg artifacts, signature, watermark, username, blurry"

image = pipe(
    prompt, 
    negative_prompt=negative_prompt, 
    width=1024,
    height=1024,
    guidance_scale=12,
    target_size=(1024,1024),
    original_size=(4096,4096),
    num_inference_steps=50
    ).images[0]

image.save("anime_girl.png")
```
<hr>

## Limitation 
This model inherit Stable Diffusion XL 1.0 [limitation](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0#limitations)