--- license: openrail++ language: - en tags: - text-to-image - stable-diffusion - lora - safetensors - diffusers - stable-diffusion-xl base_model: Linaqruf/animagine-xl inference: parameter: negative_prompt: lowres, bad anatomy, bad hands, text, error, missing fingers, extra digit, fewer digits, cropped, worst quality, low quality, normal quality, jpeg artifacts, signature, watermark, username, blurry widget: - text: >- face focus, cute, masterpiece, best quality, 1girl, green hair, sweater, looking at viewer, upper body, beanie, outdoors, night, turtleneck example_title: example 1girl - text: >- face focus, bishounen, masterpiece, best quality, 1boy, green hair, sweater, looking at viewer, upper body, beanie, outdoors, night, turtleneck example_title: example 1boy datasets: - Linaqruf/sdxl-dataset ---

Pastel Anime LoRA for SDXL

TRAINED WITH ANIMAGINE XL


sample1 sample2

## Overview **Pastel Anime LoRA for SDXL** is a high-resolution, Low-Rank Adaptation model for Stable Diffusion XL. The model has been fine-tuned using a learning rate of 1e-5 over 1300 global steps with a batch size of 24 on a curated dataset of superior-quality anime-style images. This model is derived from Animagine XL. Like other anime-style Stable Diffusion models, it also supports Danbooru tags to generate images. e.g. _**face focus, cute, masterpiece, best quality, 1girl, green hair, sweater, looking at viewer, upper body, beanie, outdoors, night, turtleneck**_
## Model Details - **Developed by:** [Linaqruf](https://github.com/Linaqruf) - **Model type:** Low-rank adaptation of diffusion-based text-to-image generative model - **Model Description:** This is a small model that should be used with big model and can be used to generate and modify high quality anime-themed images based on text prompts. - **License:** [CreativeML Open RAIL++-M License](https://huggingface.co/stabilityai/stable-diffusion-2/blob/main/LICENSE-MODEL) - **Finetuned from model:** [Animagine XL](https://huggingface.co/Linaqruf/animagine-xl)
## 🧨 Diffusers Make sure to upgrade diffusers to >= 0.18.2: ``` pip install diffusers --upgrade ``` In addition make sure to install `transformers`, `safetensors`, `accelerate` as well as the invisible watermark: ``` pip install invisible_watermark transformers accelerate safetensors ``` Running the pipeline (if you don't swap the scheduler it will run with the default **EulerDiscreteScheduler** in this example we are swapping it to **EulerAncestralDiscreteScheduler**: ```py import torch from torch import autocast from diffusers import StableDiffusionXLPipeline, EulerAncestralDiscreteScheduler base_model = "Linaqruf/animagine-xl" lora_model_id = "Linaqruf/pastel-anime-xl-lora" lora_filename = "pastel-anime-xl.safetensors" pipe = StableDiffusionXLPipeline.from_pretrained( model, torch_dtype=torch.float16, use_safetensors=True, variant="fp16" ) pipe.scheduler = EulerAncestralDiscreteScheduler.from_config(pipe.scheduler.config) pipe.to('cuda') pipe.load_lora_weights(lora_model_id, weight_name=lora_filename) prompt = "face focus, cute, masterpiece, best quality, 1girl, green hair, sweater, looking at viewer, upper body, beanie, outdoors, night, turtleneck" negative_prompt = "lowres, bad anatomy, bad hands, text, error, missing fingers, extra digit, fewer digits, cropped, worst quality, low quality, normal quality, jpeg artifacts, signature, watermark, username, blurry" image = pipe( prompt, negative_prompt=negative_prompt, width=1024, height=1024, guidance_scale=12, target_size=(1024,1024), original_size=(4096,4096), num_inference_steps=50 ).images[0] image.save("anime_girl.png") ```
## Limitation This model inherit Stable Diffusion XL 1.0 [limitation](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0#limitations)