import gradio as gr from base64 import b64encode import numpy import torch from diffusers import AutoencoderKL, LMSDiscreteScheduler, UNet2DConditionModel, StableDiffusionPipeline # For video display: from matplotlib import pyplot as plt from pathlib import Path from PIL import Image from torch import autocast from torchvision import transforms as tfms from tqdm.auto import tqdm from transformers import CLIPTextModel, CLIPTokenizer, logging import os torch.manual_seed(1) # Supress some unnecessary warnings when loading the CLIPTextModel logging.set_verbosity_error() # Set device torch_device = "cuda" if torch.cuda.is_available() else "mps" if torch.backends.mps.is_available() else "cpu" if "mps" == torch_device: os.environ['PYTORCH_ENABLE_MPS_FALLBACK'] = "1" model_nm = "CompVis/stable-diffusion-v1-4" output_dir="sd-concept-output" pipe = StableDiffusionPipeline.from_pretrained(model_nm).to(torch_device) # Load the autoencoder model which will be used to decode the latents into image space. vae = pipe.vae tokenizer = pipe.tokenizer # Load the tokenizer and text encoder to tokenize and encode the text. #tokenizer = CLIPTokenizer.from_pretrained("openai/clip-vit-large-patch14", torch_dtype=torch.float16) text_encoder =pipe.text_encoder # The UNet model for generating the latents. unet = pipe.unet # The noise scheduler scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000) # To the GPU we go! vae = vae.to(torch_device) text_encoder = text_encoder.to(torch_device) unet = unet.to(torch_device); pipe.load_textual_inversion("sd-concepts-library/madhubani-art") pipe.load_textual_inversion("sd-concepts-library/line-art") pipe.load_textual_inversion("sd-concepts-library/cat-toy") pipe.load_textual_inversion("sd-concepts-library/concept-art") def pil_to_latent(input_im): # Single image -> single latent in a batch (so size 1, 4, 64, 64) with torch.no_grad(): latent = vae.encode(tfms.ToTensor()(input_im).unsqueeze(0).to(torch_device)*2-1) # Note scaling return 0.18215 * latent.latent_dist.sample() def latents_to_pil(latents): # bath of latents -> list of images latents = (1 / 0.18215) * latents with torch.no_grad(): image = vae.decode(latents).sample image = (image / 2 + 0.5).clamp(0, 1) image = image.detach().cpu().permute(0, 2, 3, 1).numpy() images = (image * 255).round().astype("uint8") pil_images = [Image.fromarray(image) for image in images] return pil_images # Prep Scheduler def set_timesteps(scheduler, num_inference_steps): scheduler.set_timesteps(num_inference_steps) scheduler.timesteps = scheduler.timesteps.to(torch.float32) # minor fix to ensure MPS compatibility, fixed in diffusers def saturation_loss(images): # Calculate saturation for each pixel in the input image tensor max_vals, _ = torch.max(images, dim=1, keepdim=True) min_vals, _ = torch.min(images, dim=1, keepdim=True) saturation = (max_vals - min_vals) / max_vals.clamp(min=1e-7) # Avoid division by zero # Calculate mean saturation across the image mean_saturation = torch.mean(saturation, dim=(2, 3)) # Average over width and height # Calculate the loss as the negative mean saturation (proportional to saturation) #loss = torch.abs(saturation - 0.9).mean() return mean_saturation/10000 def generateImage(prompt, lossScale): #prompt = 'a puppy in style' #@param height = 512 # default height of Stable Diffusion width = 512 # default width of Stable Diffusion num_inference_steps = 200 #@param # Number of denoising steps guidance_scale = 8 #@param # Scale for classifier-free guidance generator = torch.manual_seed(32) # Seed generator to create the inital latent noise batch_size = 1 saturation_loss_Scale = lossScale #@param # Prep text text_input = tokenizer([prompt], padding="max_length", max_length=tokenizer.model_max_length, truncation=True, return_tensors="pt") with torch.no_grad(): text_embeddings = text_encoder(text_input.input_ids.to(torch_device))[0] # And the uncond. input as before: max_length = text_input.input_ids.shape[-1] uncond_input = tokenizer( [""] * batch_size, padding="max_length", max_length=max_length, return_tensors="pt" ) with torch.no_grad(): uncond_embeddings = text_encoder(uncond_input.input_ids.to(torch_device))[0] text_embeddings = torch.cat([uncond_embeddings, text_embeddings]) # Prep Scheduler set_timesteps(scheduler, num_inference_steps) # Prep latents latents = torch.randn( (batch_size, unet.in_channels, height // 8, width // 8), generator=generator, ) latents = latents.to(torch_device) latents = latents * scheduler.init_noise_sigma # Loop for i, t in tqdm(enumerate(scheduler.timesteps), total=len(scheduler.timesteps)): # expand the latents if we are doing classifier-free guidance to avoid doing two forward passes. latent_model_input = torch.cat([latents] * 2) sigma = scheduler.sigmas[i] latent_model_input = scheduler.scale_model_input(latent_model_input, t) # predict the noise residual with torch.no_grad(): noise_pred = unet(latent_model_input, t, encoder_hidden_states=text_embeddings)["sample"] # perform CFG noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) #### ADDITIONAL GUIDANCE ### if i%5 == 0: # Requires grad on the latents latents = latents.detach().requires_grad_() # Get the predicted x0: # latents_x0 = latents - sigma * noise_pred latents_x0 = scheduler.step(noise_pred, t, latents).pred_original_sample scheduler._step_index = scheduler._step_index - 1 # Decode to image space denoised_images = vae.decode((1 / 0.18215) * latents_x0).sample / 2 + 0.5 # range (0, 1) # Calculate loss loss = saturation_loss(denoised_images) * saturation_loss_Scale #loss = loss.detach().requires_grad_() #print('loss.grad_fn = {}'.format(grad_fn)) # Occasionally print it out if i%10==0: print(i, 'loss:', loss.item()) # Get gradient cond_grad = torch.autograd.grad(loss, latents)[0] # Modify the latents based on this gradient latents = latents.detach() - cond_grad * sigma**2 # Now step with scheduler latents = scheduler.step(noise_pred, t, latents).prev_sample custom_loss_image = latents_to_pil(latents)[0] return custom_loss_image def inference(imgText, style, customLoss="no"): prompt = f'a {imgText} in <{style}> style' if (customLoss == "yes") : outImage = generateImage(prompt, 2) return outImage else: outImage = generateImage(prompt, 0) return outImage title = "TSAI S20 Assignment: Use a pretrained Sstable Diffusion model and give a demo on its workig" description = "A simple Gradio interface that accepts a text and style, and generated an image using stable diffusion pipeline" examples = [["puppy","cat-toy","yes"]] demo = gr.Interface( inference, inputs = [gr.Textbox("Enter an image you want to generate"), gr.Dropdown(["madhubani-art", "line-art", "cat-toy","concept-art"], label="Choose your style"), gr.Radio(["yes", "no"], label="Add custom saturation loss?") ], outputs = [gr.Image(shape=(512, 512), label="Generated Image")], title = title, description = description, examples = examples, ) demo.launch()