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import gradio as gr | |
import os | |
import sys | |
from base64 import b64encode | |
import numpy | |
import numpy as np | |
import torch | |
from diffusers import AutoencoderKL, LMSDiscreteScheduler, UNet2DConditionModel | |
from matplotlib import pyplot as plt | |
from pathlib import Path | |
from PIL import Image | |
from torch import autocast | |
from torchvision import transforms as tfms | |
from tqdm.auto import tqdm | |
from transformers import CLIPTextModel, CLIPTokenizer, logging | |
import os | |
import cv2 | |
import torchvision.transforms as T | |
torch.manual_seed(1) | |
logging.set_verbosity_error() | |
torch_device = "cuda" if torch.cuda.is_available() else "cpu" | |
# Load the autoencoder | |
vae = AutoencoderKL.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder='vae') | |
# Load tokenizer and text encoder to tokenize and encode the text | |
tokenizer = CLIPTokenizer.from_pretrained("openai/clip-vit-large-patch14") | |
text_encoder = CLIPTextModel.from_pretrained("openai/clip-vit-large-patch14") | |
# Unet model for generating latents | |
unet = UNet2DConditionModel.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder='unet') | |
# Noise scheduler | |
scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000) | |
# Move everything to GPU | |
vae = vae.to(torch_device) | |
text_encoder = text_encoder.to(torch_device) | |
unet = unet.to(torch_device) | |
def get_output_embeds(input_embeddings): | |
# CLIP's text model uses causal mask, so we prepare it here: | |
bsz, seq_len = input_embeddings.shape[:2] | |
causal_attention_mask = text_encoder.text_model._build_causal_attention_mask(bsz, seq_len, dtype=input_embeddings.dtype) | |
# Getting the output embeddings involves calling the model with passing output_hidden_states=True | |
# so that it doesn't just return the pooled final predictions: | |
encoder_outputs = text_encoder.text_model.encoder( | |
inputs_embeds=input_embeddings, | |
attention_mask=None, # We aren't using an attention mask so that can be None | |
causal_attention_mask=causal_attention_mask.to(torch_device), | |
output_attentions=None, | |
output_hidden_states=True, # We want the output embs not the final output | |
return_dict=None, | |
) | |
# We're interested in the output hidden state only | |
output = encoder_outputs[0] | |
# There is a final layer norm we need to pass these through | |
output = text_encoder.text_model.final_layer_norm(output) | |
# And now they're ready! | |
return output | |
# Prep Scheduler | |
def set_timesteps(scheduler, num_inference_steps): | |
scheduler.set_timesteps(num_inference_steps) | |
scheduler.timesteps = scheduler.timesteps.to(torch.float32) # minor fix to ensure MPS compatibility, fixed in diffusers PR 3925 | |
style_files = ['learned_embeds_birb_style.bin','learned_embeds_cute_game_style.bin', | |
'learned_embeds_manga_style.bin','learned_embeds_midjourney_style.bin','learned_embeds_space_style.bin'] | |
seed_values = [8,16,50,80,128] | |
height = 512 # default height of Stable Diffusion | |
width = 512 # default width of Stable Diffusion | |
num_inference_steps = 5 # Number of denoising steps | |
guidance_scale = 7.5 # Scale for classifier-free guidance | |
num_styles = len(style_files) | |
def get_style_embeddings(style_file): | |
style_embed = torch.load(style_file) | |
style_name = list(style_embed.keys())[0] | |
return style_embed[style_name] | |
def get_EOS_pos_in_prompt(prompt): | |
return len(prompt.split())+1 | |
import torch.nn.functional as F | |
""" | |
def gradient_loss(images): | |
# Compute gradient magnitude using Sobel filters. | |
gradient_x = F.conv2d(images, torch.Tensor([[-1, 0, 1], [-2, 0, 2], [-1, 0, 1]]).view(1, 1, 3, 3).to(images.device)) | |
gradient_y = F.conv2d(images, torch.Tensor([[-1, -2, -1], [0, 0, 0], [1, 2, 1]]).view(1, 1, 3, 3).to(images.device)) | |
gradient_magnitude = torch.sqrt(gradient_x**2 + gradient_y**2) | |
return gradient_magnitude.mean() | |
""" | |
from torchvision.transforms import ToTensor | |
def pil_to_latent(input_im): | |
# Single image -> single latent in a batch (so size 1, 4, 64, 64) | |
with torch.no_grad(): | |
latent = vae.encode(tfms.ToTensor()(input_im).unsqueeze(0).to(torch_device)*2-1) # Note scaling | |
return 0.18215 * latent.latent_dist.sample() | |
def latents_to_pil(latents): | |
# bath of latents -> list of images | |
latents = (1 / 0.18215) * latents | |
with torch.no_grad(): | |
image = vae.decode(latents).sample | |
image = (image / 2 + 0.5).clamp(0, 1) | |
image = image.detach().cpu().permute(0, 2, 3, 1).numpy() | |
images = (image * 255).round().astype("uint8") | |
pil_images = [Image.fromarray(image) for image in images] | |
return pil_images | |
def additional_guidance(latents, scheduler, noise_pred, t, sigma, custom_loss_fn, custom_loss_scale): | |
#### ADDITIONAL GUIDANCE ### | |
# Requires grad on the latents | |
latents = latents.detach().requires_grad_() | |
# Get the predicted x0: | |
latents_x0 = latents - sigma * noise_pred | |
# Decode to image space | |
denoised_images = vae.decode((1 / 0.18215) * latents_x0).sample / 2 + 0.5 # range (0, 1) | |
# Calculate loss | |
loss = custom_loss_fn(denoised_images) * custom_loss_scale | |
# Get gradient | |
cond_grad = torch.autograd.grad(loss, latents, allow_unused=False)[0] | |
# Modify the latents based on this gradient | |
latents = latents.detach() - cond_grad * sigma**2 | |
return latents, loss | |
def generate_with_embs(text_embeddings, max_length, random_seed, loss_fn = None, custom_loss_scale=1.0): | |
height = 512 # default height of Stable Diffusion | |
width = 512 # default width of Stable Diffusion | |
num_inference_steps = 5 # Number of denoising steps | |
guidance_scale = 7.5 # Scale for classifier-free guidance | |
generator = torch.manual_seed(random_seed) # Seed generator to create the inital latent noise | |
batch_size = 1 | |
uncond_input = tokenizer( | |
[""] * batch_size, padding="max_length", max_length=max_length, return_tensors="pt" | |
) | |
with torch.no_grad(): | |
uncond_embeddings = text_encoder(uncond_input.input_ids.to(torch_device))[0] | |
text_embeddings = torch.cat([uncond_embeddings, text_embeddings]) | |
# Prep Scheduler | |
set_timesteps(scheduler, num_inference_steps) | |
# Prep latents | |
latents = torch.randn( | |
(batch_size, unet.in_channels, height // 8, width // 8), | |
generator=generator, | |
) | |
latents = latents.to(torch_device) | |
latents = latents * scheduler.init_noise_sigma | |
# Loop | |
for i, t in tqdm(enumerate(scheduler.timesteps), total=len(scheduler.timesteps)): | |
# expand the latents if we are doing classifier-free guidance to avoid doing two forward passes. | |
latent_model_input = torch.cat([latents] * 2) | |
sigma = scheduler.sigmas[i] | |
latent_model_input = scheduler.scale_model_input(latent_model_input, t) | |
# predict the noise residual | |
with torch.no_grad(): | |
noise_pred = unet(latent_model_input, t, encoder_hidden_states=text_embeddings)["sample"] | |
# perform guidance | |
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) | |
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) | |
if loss_fn is not None: | |
if i%2 == 0: | |
latents, custom_loss = additional_guidance(latents, scheduler, noise_pred, t, sigma, loss_fn, custom_loss_scale) | |
print(i, 'loss:', custom_loss.item()) | |
# compute the previous noisy sample x_t -> x_t-1 | |
latents = scheduler.step(noise_pred, t, latents).prev_sample | |
return latents_to_pil(latents)[0] | |
def generate_image_custom_style(prompt, style_num=None, random_seed=41, custom_loss_fn = None, custom_loss_scale=1.0): | |
eos_pos = get_EOS_pos_in_prompt(prompt) | |
style_token_embedding = None | |
if style_num: | |
style_token_embedding = get_style_embeddings(style_files[style_num]) | |
# tokenize | |
text_input = tokenizer(prompt, padding="max_length", max_length=tokenizer.model_max_length, truncation=True, return_tensors="pt") | |
max_length = text_input.input_ids.shape[-1] | |
input_ids = text_input.input_ids.to(torch_device) | |
# get token embeddings | |
token_emb_layer = text_encoder.text_model.embeddings.token_embedding | |
token_embeddings = token_emb_layer(input_ids) | |
# Append style token towards the end of the sentence embeddings | |
if style_token_embedding is not None: | |
token_embeddings[-1, eos_pos, :] = style_token_embedding | |
# combine with pos embs | |
pos_emb_layer = text_encoder.text_model.embeddings.position_embedding | |
position_ids = text_encoder.text_model.embeddings.position_ids[:, :77] | |
position_embeddings = pos_emb_layer(position_ids) | |
input_embeddings = token_embeddings + position_embeddings | |
# Feed through to get final output embs | |
modified_output_embeddings = get_output_embeds(input_embeddings) | |
# And generate an image with this: | |
generated_image = generate_with_embs(modified_output_embeddings, max_length, random_seed, custom_loss_fn, custom_loss_scale) | |
return generated_image | |
def show_images(images_list): | |
# Let's visualize the four channels of this latent representation: | |
fig, axs = plt.subplots(1, len(images_list), figsize=(16, 4)) | |
for c in range(len(images_list)): | |
axs[c].imshow(images_list[c]) | |
plt.show() | |
def invert_loss(gen_image): | |
inverter = T.RandomInvert(p=1.0) | |
inverted_img = inverter(gen_image) | |
#loss = torch.abs(gen_image - inverted_img).sum() | |
loss = torch.nn.functional.mse_loss(gen_image[:,0], gen_image[:,2]) + torch.nn.functional.mse_loss(gen_image[:,2], gen_image[:,1]) + torch.nn.functional.mse_loss(gen_image[:,0], gen_image[:,1]) | |
return loss | |
def contrast_loss(images): | |
# Calculate the variance of pixel values as a measure of contrast. | |
variance = torch.var(images) | |
return -variance | |
def blue_loss(images): | |
# How far are the blue channel values to 0.9: | |
error = torch.abs(images[:,2] - 0.9).mean() # [:,2] -> all images in batch, only the blue channel | |
return error | |
def display_images_in_rows(images_with_titles, titles): | |
num_images = len(images_with_titles) | |
rows = 5 # Display 5 rows always | |
columns = 1 if num_images == 5 else 2 # Use 1 column if there are 5 images, otherwise 2 columns | |
fig, axes = plt.subplots(rows, columns + 1, figsize=(15, 5 * rows)) # Add an extra column for titles | |
for r in range(rows): | |
# Add the title on the extreme left in the middle of each picture | |
axes[r, 0].text(0.5, 0.5, titles[r], ha='center', va='center') | |
axes[r, 0].axis('off') | |
# Add "Without Loss" label above the first column and "With Loss" label above the second column (if applicable) | |
if columns == 2: | |
axes[r, 1].set_title("Without Loss", pad=10) | |
axes[r, 2].set_title("With Loss", pad=10) | |
for c in range(1, columns + 1): | |
index = r * columns + c - 1 | |
if index < num_images: | |
image, _ = images_with_titles[index] | |
axes[r, c].imshow(image) | |
axes[r, c].axis('off') | |
return fig | |
# plt.show() | |
def image_generator(prompt = "dog", loss_function=None): | |
images_without_loss = [] | |
images_with_loss = [] | |
for i in range(num_styles): | |
generated_img = generate_image_custom_style(prompt,style_num = i,random_seed = seed_values[i],custom_loss_fn = None) | |
images_without_loss.append(generated_img) | |
if loss_function: | |
generated_img = generate_image_custom_style(prompt,style_num = i,random_seed = seed_values[i],custom_loss_fn = loss_function) | |
images_with_loss.append(generated_img) | |
generated_sd_images = [] | |
titles = ["Birb Style","Cute Game Style","Manga Style","Mid Journey Style","Space Style"] | |
for i in range(len(titles)): | |
generated_sd_images.append((images_without_loss[i], titles[i])) | |
if images_with_loss != []: | |
generated_sd_images.append((images_with_loss[i], titles[i])) | |
return display_images_in_rows(generated_sd_images, titles) | |
# Create a wrapper function for show_misclassified_images() | |
def image_generator_wrapper(prompt = "dog", loss_function=None): | |
if loss_function == "Yes": | |
loss_function = contrast_loss | |
else: | |
loss_function = None | |
return image_generator(prompt, loss_function) | |
description = '(Team Project EE267) Stable Diffusion is a generative artificial intelligence (generative AI) model that produces unique photorealistic images from text and image prompts.' | |
title = 'Image Generation using Stable Diffusion' | |
demo = gr.Interface(image_generator_wrapper, | |
inputs=[gr.Textbox(label="Enter prompt for generation", type="text", value="a dog walking on moon"), | |
gr.Radio(["Yes", "No"], value="No" , label="Custom Loss Function")], | |
outputs=gr.Plot(label="Generated Images"), title = "EE 267 Team Project - Stable Diffusion", description=description) | |
demo.launch() |