license: openrail++
language:
- en
tags:
- stable-diffusion
- text-to-image
- diffusers
thumbnail: >-
https://huggingface.co/nitrosocke/Future-Diffusion/resolve/main/images/future-diffusion-thumbnail-2.jpg
inference: false
Future Diffusion
This is the fine-tuned Stable Diffusion 2.0 model trained on high quality 3D images with a 768x768 pixel resolution.
Use the tokensredshift style
in your prompts for the effect.
Trained on Stability.ai's Stable Diffusion 2.0 with 768x768 resolution.
If you enjoy my work and want to test new models before release, please consider supporting me
Disclaimer: The SD 2.0 model is just a few days old at this point and we still need to figure out how it works exactly. Please view this as an early prototype and experiment with the model.
Characters rendered with the model:
Cars and Animals rendered with the model:
Landscapes rendered with the model:
Prompt and settings for the Characters:
redshift style [subject] Negative Prompt:
Steps: 20, Sampler: Euler a, CFG scale: 7, Size: 768x1024
Prompt and settings for the Landscapes:
redshift style Negative Prompt:
Steps: 20, Sampler: Euler a, CFG scale: 7, Size: 1152x768
This model was trained using the diffusers based dreambooth training by ShivamShrirao using prior-preservation loss and the train-text-encoder flag in 7.500 steps.
License
This model is open access and available to all, with a CreativeML Open RAIL++-M License further specifying rights and usage. Please read the full license here